7.8 KiB
Stable Diffusion
Stable Diffusion is an open-source diffusion-based text-to-image generation model developed by Stability AI, supporting 512x512 resolution text-to-image generation.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .
For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will quickly load the AI-ModelScope/stable-diffusion-v1-5 model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 2GB VRAM.
import torch
from diffsynth.core import ModelConfig
from diffsynth.pipelines.stable_diffusion import StableDiffusionPipeline
vram_config = {
"offload_dtype": torch.float32,
"offload_device": "cpu",
"onload_dtype": torch.float32,
"onload_device": "cpu",
"preparing_dtype": torch.float32,
"preparing_device": "cuda",
"computation_dtype": torch.float32,
"computation_device": "cuda",
}
pipe = StableDiffusionPipeline.from_pretrained(
torch_dtype=torch.float32,
model_configs=[
ModelConfig(model_id="AI-ModelScope/stable-diffusion-v1-5", origin_file_pattern="text_encoder/model.safetensors", **vram_config),
ModelConfig(model_id="AI-ModelScope/stable-diffusion-v1-5", origin_file_pattern="unet/diffusion_pytorch_model.safetensors", **vram_config),
ModelConfig(model_id="AI-ModelScope/stable-diffusion-v1-5", origin_file_pattern="vae/diffusion_pytorch_model.safetensors", **vram_config),
],
tokenizer_config=ModelConfig(model_id="AI-ModelScope/stable-diffusion-v1-5", origin_file_pattern="tokenizer/"),
vram_limit=torch.cuda.mem_get_info("cuda")[1] / (1024 ** 3) - 0.5,
)
image = pipe(
prompt="a photo of an astronaut riding a horse on mars, high quality, detailed",
negative_prompt="blurry, low quality, deformed",
cfg_scale=7.5,
height=512,
width=512,
seed=42,
rand_device="cuda",
num_inference_steps=50,
)
image.save("image.jpg")
Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| AI-ModelScope/stable-diffusion-v1-5 | code | code | code | code | code | code |
Model Inference
The model is loaded via StableDiffusionPipeline.from_pretrained, see Loading Models for details.
The input parameters for StableDiffusionPipeline inference include:
prompt: Text prompt.negative_prompt: Negative prompt, defaults to an empty string.cfg_scale: Classifier-Free Guidance scale factor, default 7.5.height: Output image height, default 512.width: Output image width, default 512.seed: Random seed, defaults to a random value if not set.rand_device: Noise generation device, defaults to "cpu".num_inference_steps: Number of inference steps, default 50.eta: DDIM scheduler eta parameter, default 0.0.guidance_rescale: Guidance rescale factor, default 0.0.progress_bar_cmd: Progress bar callback function.
Model Training
Models in the stable_diffusion series are trained via examples/stable_diffusion/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- Stable Diffusion Specific Parameters
--tokenizer_path: Tokenizer path, defaults toAI-ModelScope/stable-diffusion-v1-5:tokenizer/.
Example dataset download:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --include "stable_diffusion/*" --local_dir ./data/diffsynth_example_dataset
stable-diffusion-v1-5 training scripts
We provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.